Stable Diffusion v1-5
Author runwayml
model_id: sd-1.5
Model Type: Stable Diffusion
Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. For more information about h...
API Code Snippet
curl --location --request POST 'https://stablediffusionapi.com/api/v4/dreambooth' \
curl --header 'Content-Type: application/json' \
{
"key": "api-key",
"model_id": "sd-1.5",
"prompt": "ultra realistic close up portrait ((beautiful pale cyberpunk female with heavy black eyeliner)), blue eyes, shaved side haircut, hyper detail, cinematic lighting, magic neon, dark red city, Canon EOS R3, nikon, f/1.4, ISO 200, 1/160s, 8K, RAW, unedited, symmetrical balance, in-frame, 8K",
"negative_prompt": "painting, extra fingers, mutated hands, poorly drawn hands, poorly drawn face, deformed, ugly, blurry, bad anatomy, bad proportions, extra limbs, cloned face, skinny, glitchy, double torso, extra arms, extra hands, mangled fingers, missing lips, ugly face, distorted face, extra legs, anime",
"width": "512",
"height": "512",
"samples": "1",
"num_inference_steps": "30",
"seed": null,
"guidance_scale": 7.5,
"webhook": null,
"track_id": null
}
No. of images generated
12.01 K
Images
Generated
Images generated with Stable Diffusion v1-5 and its prompt